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| # Copyright 2023 The HuggingFace Team. All rights reserved. | |
| # | |
| # Licensed under the Apache License, Version 2.0 (the "License"); | |
| # you may not use this file except in compliance with the License. | |
| # You may obtain a copy of the License at | |
| # | |
| # http://www.apache.org/licenses/LICENSE-2.0 | |
| # | |
| # Unless required by applicable law or agreed to in writing, software | |
| # distributed under the License is distributed on an "AS IS" BASIS, | |
| # WITHOUT WARRANTIES OR CONDITIONS OF ANY KIND, either express or implied. | |
| # See the License for the specific language governing permissions and | |
| # limitations under the License. | |
| import inspect | |
| from typing import Any, Callable, Dict, List, Optional, Tuple, Union | |
| import numpy as np | |
| import PIL.Image | |
| import torch | |
| from transformers import CLIPTextModel, CLIPTextModelWithProjection, CLIPTokenizer | |
| import torchvision | |
| from diffusers.image_processor import VaeImageProcessor | |
| from diffusers.loaders import FromSingleFileMixin, LoraLoaderMixin, TextualInversionLoaderMixin | |
| from diffusers.models import AutoencoderKL, UNet2DConditionModel | |
| from diffusers.models.attention_processor import ( | |
| AttnProcessor2_0, | |
| LoRAAttnProcessor2_0, | |
| LoRAXFormersAttnProcessor, | |
| XFormersAttnProcessor, | |
| ) | |
| from diffusers.schedulers import KarrasDiffusionSchedulers | |
| from diffusers.utils import ( | |
| is_accelerate_available, | |
| is_accelerate_version, | |
| is_invisible_watermark_available, | |
| logging, | |
| randn_tensor, | |
| replace_example_docstring, | |
| ) | |
| from diffusers.pipelines.pipeline_utils import DiffusionPipeline | |
| from diffusers.pipelines.stable_diffusion_xl import StableDiffusionXLPipelineOutput | |
| if is_invisible_watermark_available(): | |
| from .watermark import StableDiffusionXLWatermarker | |
| logger = logging.get_logger(__name__) # pylint: disable=invalid-name | |
| EXAMPLE_DOC_STRING = """ | |
| Examples: | |
| ```py | |
| >>> import torch | |
| >>> from diffusers import StableDiffusionXLImg2ImgPipeline | |
| >>> from diffusers.utils import load_image | |
| >>> pipe = StableDiffusionXLImg2ImgPipeline.from_pretrained( | |
| ... "stabilityai/stable-diffusion-xl-refiner-1.0", torch_dtype=torch.float16 | |
| ... ) | |
| >>> pipe = pipe.to("cuda") | |
| >>> url = "https://huggingface.co/datasets/patrickvonplaten/images/resolve/main/aa_xl/000000009.png" | |
| >>> init_image = load_image(url).convert("RGB") | |
| >>> prompt = "a photo of an astronaut riding a horse on mars" | |
| >>> image = pipe(prompt, image=init_image).images[0] | |
| ``` | |
| """ | |
| # Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.rescale_noise_cfg | |
| def rescale_noise_cfg(noise_cfg, noise_pred_text, guidance_rescale=0.0): | |
| """ | |
| Rescale `noise_cfg` according to `guidance_rescale`. Based on findings of [Common Diffusion Noise Schedules and | |
| Sample Steps are Flawed](https://arxiv.org/pdf/2305.08891.pdf). See Section 3.4 | |
| """ | |
| std_text = noise_pred_text.std(dim=list(range(1, noise_pred_text.ndim)), keepdim=True) | |
| std_cfg = noise_cfg.std(dim=list(range(1, noise_cfg.ndim)), keepdim=True) | |
| # rescale the results from guidance (fixes overexposure) | |
| noise_pred_rescaled = noise_cfg * (std_text / std_cfg) | |
| # mix with the original results from guidance by factor guidance_rescale to avoid "plain looking" images | |
| noise_cfg = guidance_rescale * noise_pred_rescaled + (1 - guidance_rescale) * noise_cfg | |
| return noise_cfg | |
| class StableDiffusionXLDiffImg2ImgPipeline(DiffusionPipeline, FromSingleFileMixin, LoraLoaderMixin): | |
| r""" | |
| Pipeline for text-to-image generation using Stable Diffusion XL. | |
| This model inherits from [`DiffusionPipeline`]. Check the superclass documentation for the generic methods the | |
| library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.) | |
| In addition the pipeline inherits the following loading methods: | |
| - *Textual-Inversion*: [`loaders.TextualInversionLoaderMixin.load_textual_inversion`] | |
| - *LoRA*: [`loaders.LoraLoaderMixin.load_lora_weights`] | |
| - *Ckpt*: [`loaders.FromSingleFileMixin.from_single_file`] | |
| as well as the following saving methods: | |
| - *LoRA*: [`loaders.LoraLoaderMixin.save_lora_weights`] | |
| Args: | |
| vae ([`AutoencoderKL`]): | |
| Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations. | |
| text_encoder ([`CLIPTextModel`]): | |
| Frozen text-encoder. Stable Diffusion XL uses the text portion of | |
| [CLIP](https://huggingface.co/docs/transformers/model_doc/clip#transformers.CLIPTextModel), specifically | |
| the [clip-vit-large-patch14](https://huggingface.co/openai/clip-vit-large-patch14) variant. | |
| text_encoder_2 ([` CLIPTextModelWithProjection`]): | |
| Second frozen text-encoder. Stable Diffusion XL uses the text and pool portion of | |
| [CLIP](https://huggingface.co/docs/transformers/model_doc/clip#transformers.CLIPTextModelWithProjection), | |
| specifically the | |
| [laion/CLIP-ViT-bigG-14-laion2B-39B-b160k](https://huggingface.co/laion/CLIP-ViT-bigG-14-laion2B-39B-b160k) | |
| variant. | |
| tokenizer (`CLIPTokenizer`): | |
| Tokenizer of class | |
| [CLIPTokenizer](https://huggingface.co/docs/transformers/v4.21.0/en/model_doc/clip#transformers.CLIPTokenizer). | |
| tokenizer_2 (`CLIPTokenizer`): | |
| Second Tokenizer of class | |
| [CLIPTokenizer](https://huggingface.co/docs/transformers/v4.21.0/en/model_doc/clip#transformers.CLIPTokenizer). | |
| unet ([`UNet2DConditionModel`]): Conditional U-Net architecture to denoise the encoded image latents. | |
| scheduler ([`SchedulerMixin`]): | |
| A scheduler to be used in combination with `unet` to denoise the encoded image latents. Can be one of | |
| [`DDIMScheduler`], [`LMSDiscreteScheduler`], or [`PNDMScheduler`]. | |
| """ | |
| _optional_components = ["tokenizer", "text_encoder"] | |
| def __init__( | |
| self, | |
| vae: AutoencoderKL, | |
| text_encoder: CLIPTextModel, | |
| text_encoder_2: CLIPTextModelWithProjection, | |
| tokenizer: CLIPTokenizer, | |
| tokenizer_2: CLIPTokenizer, | |
| unet: UNet2DConditionModel, | |
| scheduler: KarrasDiffusionSchedulers, | |
| requires_aesthetics_score: bool = False, | |
| force_zeros_for_empty_prompt: bool = True, | |
| add_watermarker: Optional[bool] = None, | |
| ): | |
| super().__init__() | |
| self.register_modules( | |
| vae=vae, | |
| text_encoder=text_encoder, | |
| text_encoder_2=text_encoder_2, | |
| tokenizer=tokenizer, | |
| tokenizer_2=tokenizer_2, | |
| unet=unet, | |
| scheduler=scheduler, | |
| ) | |
| self.register_to_config(force_zeros_for_empty_prompt=force_zeros_for_empty_prompt) | |
| self.register_to_config(requires_aesthetics_score=requires_aesthetics_score) | |
| self.vae_scale_factor = 2 ** (len(self.vae.config.block_out_channels) - 1) | |
| self.image_processor = VaeImageProcessor(vae_scale_factor=self.vae_scale_factor) | |
| add_watermarker = add_watermarker if add_watermarker is not None else is_invisible_watermark_available() | |
| if add_watermarker: | |
| self.watermark = StableDiffusionXLWatermarker() | |
| else: | |
| self.watermark = None | |
| # Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_vae_slicing | |
| def enable_vae_slicing(self): | |
| r""" | |
| Enable sliced VAE decoding. When this option is enabled, the VAE will split the input tensor in slices to | |
| compute decoding in several steps. This is useful to save some memory and allow larger batch sizes. | |
| """ | |
| self.vae.enable_slicing() | |
| # Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_vae_slicing | |
| def disable_vae_slicing(self): | |
| r""" | |
| Disable sliced VAE decoding. If `enable_vae_slicing` was previously enabled, this method will go back to | |
| computing decoding in one step. | |
| """ | |
| self.vae.disable_slicing() | |
| # Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.enable_vae_tiling | |
| def enable_vae_tiling(self): | |
| r""" | |
| Enable tiled VAE decoding. When this option is enabled, the VAE will split the input tensor into tiles to | |
| compute decoding and encoding in several steps. This is useful for saving a large amount of memory and to allow | |
| processing larger images. | |
| """ | |
| self.vae.enable_tiling() | |
| # Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.disable_vae_tiling | |
| def disable_vae_tiling(self): | |
| r""" | |
| Disable tiled VAE decoding. If `enable_vae_tiling` was previously enabled, this method will go back to | |
| computing decoding in one step. | |
| """ | |
| self.vae.disable_tiling() | |
| def enable_model_cpu_offload(self, gpu_id=0): | |
| r""" | |
| Offloads all models to CPU using accelerate, reducing memory usage with a low impact on performance. Compared | |
| to `enable_sequential_cpu_offload`, this method moves one whole model at a time to the GPU when its `forward` | |
| method is called, and the model remains in GPU until the next model runs. Memory savings are lower than with | |
| `enable_sequential_cpu_offload`, but performance is much better due to the iterative execution of the `unet`. | |
| """ | |
| if is_accelerate_available() and is_accelerate_version(">=", "0.17.0.dev0"): | |
| from accelerate import cpu_offload_with_hook | |
| else: | |
| raise ImportError("`enable_model_cpu_offload` requires `accelerate v0.17.0` or higher.") | |
| device = torch.device(f"cuda:{gpu_id}") | |
| if self.device.type != "cpu": | |
| self.to("cpu", silence_dtype_warnings=True) | |
| torch.cuda.empty_cache() # otherwise we don't see the memory savings (but they probably exist) | |
| model_sequence = ( | |
| [self.text_encoder, self.text_encoder_2] if self.text_encoder is not None else [self.text_encoder_2] | |
| ) | |
| model_sequence.extend([self.unet, self.vae]) | |
| hook = None | |
| for cpu_offloaded_model in model_sequence: | |
| _, hook = cpu_offload_with_hook(cpu_offloaded_model, device, prev_module_hook=hook) | |
| # We'll offload the last model manually. | |
| self.final_offload_hook = hook | |
| # Copied from diffusers.pipelines.stable_diffusion_xl.pipeline_stable_diffusion_xl.StableDiffusionXLPipeline.encode_prompt | |
| def encode_prompt( | |
| self, | |
| prompt: str, | |
| prompt_2: Optional[str] = None, | |
| device: Optional[torch.device] = None, | |
| num_images_per_prompt: int = 1, | |
| do_classifier_free_guidance: bool = True, | |
| negative_prompt: Optional[str] = None, | |
| negative_prompt_2: Optional[str] = None, | |
| prompt_embeds: Optional[torch.FloatTensor] = None, | |
| negative_prompt_embeds: Optional[torch.FloatTensor] = None, | |
| pooled_prompt_embeds: Optional[torch.FloatTensor] = None, | |
| negative_pooled_prompt_embeds: Optional[torch.FloatTensor] = None, | |
| lora_scale: Optional[float] = None, | |
| ): | |
| r""" | |
| Encodes the prompt into text encoder hidden states. | |
| Args: | |
| prompt (`str` or `List[str]`, *optional*): | |
| prompt to be encoded | |
| prompt_2 (`str` or `List[str]`, *optional*): | |
| The prompt or prompts to be sent to the `tokenizer_2` and `text_encoder_2`. If not defined, `prompt` is | |
| used in both text-encoders | |
| device: (`torch.device`): | |
| torch device | |
| num_images_per_prompt (`int`): | |
| number of images that should be generated per prompt | |
| do_classifier_free_guidance (`bool`): | |
| whether to use classifier free guidance or not | |
| negative_prompt (`str` or `List[str]`, *optional*): | |
| The prompt or prompts not to guide the image generation. If not defined, one has to pass | |
| `negative_prompt_embeds` instead. Ignored when not using guidance (i.e., ignored if `guidance_scale` is | |
| less than `1`). | |
| negative_prompt_2 (`str` or `List[str]`, *optional*): | |
| The prompt or prompts not to guide the image generation to be sent to `tokenizer_2` and | |
| `text_encoder_2`. If not defined, `negative_prompt` is used in both text-encoders | |
| prompt_embeds (`torch.FloatTensor`, *optional*): | |
| Pre-generated text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. If not | |
| provided, text embeddings will be generated from `prompt` input argument. | |
| negative_prompt_embeds (`torch.FloatTensor`, *optional*): | |
| Pre-generated negative text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt | |
| weighting. If not provided, negative_prompt_embeds will be generated from `negative_prompt` input | |
| argument. | |
| pooled_prompt_embeds (`torch.FloatTensor`, *optional*): | |
| Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. | |
| If not provided, pooled text embeddings will be generated from `prompt` input argument. | |
| negative_pooled_prompt_embeds (`torch.FloatTensor`, *optional*): | |
| Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt | |
| weighting. If not provided, pooled negative_prompt_embeds will be generated from `negative_prompt` | |
| input argument. | |
| lora_scale (`float`, *optional*): | |
| A lora scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded. | |
| """ | |
| device = device or self._execution_device | |
| # set lora scale so that monkey patched LoRA | |
| # function of text encoder can correctly access it | |
| if lora_scale is not None and isinstance(self, LoraLoaderMixin): | |
| self._lora_scale = lora_scale | |
| if prompt is not None and isinstance(prompt, str): | |
| batch_size = 1 | |
| elif prompt is not None and isinstance(prompt, list): | |
| batch_size = len(prompt) | |
| else: | |
| batch_size = prompt_embeds.shape[0] | |
| # Define tokenizers and text encoders | |
| tokenizers = [self.tokenizer, self.tokenizer_2] if self.tokenizer is not None else [self.tokenizer_2] | |
| text_encoders = ( | |
| [self.text_encoder, self.text_encoder_2] if self.text_encoder is not None else [self.text_encoder_2] | |
| ) | |
| if prompt_embeds is None: | |
| prompt_2 = prompt_2 or prompt | |
| # textual inversion: procecss multi-vector tokens if necessary | |
| prompt_embeds_list = [] | |
| prompts = [prompt, prompt_2] | |
| for prompt, tokenizer, text_encoder in zip(prompts, tokenizers, text_encoders): | |
| if isinstance(self, TextualInversionLoaderMixin): | |
| prompt = self.maybe_convert_prompt(prompt, tokenizer) | |
| text_inputs = tokenizer( | |
| prompt, | |
| padding="max_length", | |
| max_length=tokenizer.model_max_length, | |
| truncation=True, | |
| return_tensors="pt", | |
| ) | |
| text_input_ids = text_inputs.input_ids | |
| untruncated_ids = tokenizer(prompt, padding="longest", return_tensors="pt").input_ids | |
| untruncated_ids = tokenizer(prompt, padding="longest", return_tensors="pt").input_ids | |
| if untruncated_ids.shape[-1] >= text_input_ids.shape[-1] and not torch.equal( | |
| text_input_ids, untruncated_ids | |
| ): | |
| removed_text = tokenizer.batch_decode(untruncated_ids[:, tokenizer.model_max_length - 1 : -1]) | |
| logger.warning( | |
| "The following part of your input was truncated because CLIP can only handle sequences up to" | |
| f" {tokenizer.model_max_length} tokens: {removed_text}" | |
| ) | |
| prompt_embeds = text_encoder( | |
| text_input_ids.to(device), | |
| output_hidden_states=True, | |
| ) | |
| # We are only ALWAYS interested in the pooled output of the final text encoder | |
| pooled_prompt_embeds = prompt_embeds[0] | |
| prompt_embeds = prompt_embeds.hidden_states[-2] | |
| prompt_embeds_list.append(prompt_embeds) | |
| prompt_embeds = torch.concat(prompt_embeds_list, dim=-1) | |
| # get unconditional embeddings for classifier free guidance | |
| zero_out_negative_prompt = negative_prompt is None and self.config.force_zeros_for_empty_prompt | |
| if do_classifier_free_guidance and negative_prompt_embeds is None and zero_out_negative_prompt: | |
| negative_prompt_embeds = torch.zeros_like(prompt_embeds) | |
| negative_pooled_prompt_embeds = torch.zeros_like(pooled_prompt_embeds) | |
| elif do_classifier_free_guidance and negative_prompt_embeds is None: | |
| negative_prompt = negative_prompt or "" | |
| negative_prompt_2 = negative_prompt_2 or negative_prompt | |
| uncond_tokens: List[str] | |
| if prompt is not None and type(prompt) is not type(negative_prompt): | |
| raise TypeError( | |
| f"`negative_prompt` should be the same type to `prompt`, but got {type(negative_prompt)} !=" | |
| f" {type(prompt)}." | |
| ) | |
| elif isinstance(negative_prompt, str): | |
| uncond_tokens = [negative_prompt, negative_prompt_2] | |
| elif batch_size != len(negative_prompt): | |
| raise ValueError( | |
| f"`negative_prompt`: {negative_prompt} has batch size {len(negative_prompt)}, but `prompt`:" | |
| f" {prompt} has batch size {batch_size}. Please make sure that passed `negative_prompt` matches" | |
| " the batch size of `prompt`." | |
| ) | |
| else: | |
| uncond_tokens = [negative_prompt, negative_prompt_2] | |
| negative_prompt_embeds_list = [] | |
| for negative_prompt, tokenizer, text_encoder in zip(uncond_tokens, tokenizers, text_encoders): | |
| if isinstance(self, TextualInversionLoaderMixin): | |
| negative_prompt = self.maybe_convert_prompt(negative_prompt, tokenizer) | |
| max_length = prompt_embeds.shape[1] | |
| uncond_input = tokenizer( | |
| negative_prompt, | |
| padding="max_length", | |
| max_length=max_length, | |
| truncation=True, | |
| return_tensors="pt", | |
| ) | |
| negative_prompt_embeds = text_encoder( | |
| uncond_input.input_ids.to(device), | |
| output_hidden_states=True, | |
| ) | |
| # We are only ALWAYS interested in the pooled output of the final text encoder | |
| negative_pooled_prompt_embeds = negative_prompt_embeds[0] | |
| negative_prompt_embeds = negative_prompt_embeds.hidden_states[-2] | |
| negative_prompt_embeds_list.append(negative_prompt_embeds) | |
| negative_prompt_embeds = torch.concat(negative_prompt_embeds_list, dim=-1) | |
| prompt_embeds = prompt_embeds.to(dtype=self.text_encoder_2.dtype, device=device) | |
| bs_embed, seq_len, _ = prompt_embeds.shape | |
| # duplicate text embeddings for each generation per prompt, using mps friendly method | |
| prompt_embeds = prompt_embeds.repeat(1, num_images_per_prompt, 1) | |
| prompt_embeds = prompt_embeds.view(bs_embed * num_images_per_prompt, seq_len, -1) | |
| if do_classifier_free_guidance: | |
| # duplicate unconditional embeddings for each generation per prompt, using mps friendly method | |
| seq_len = negative_prompt_embeds.shape[1] | |
| negative_prompt_embeds = negative_prompt_embeds.to(dtype=self.text_encoder_2.dtype, device=device) | |
| negative_prompt_embeds = negative_prompt_embeds.repeat(1, num_images_per_prompt, 1) | |
| negative_prompt_embeds = negative_prompt_embeds.view(batch_size * num_images_per_prompt, seq_len, -1) | |
| pooled_prompt_embeds = pooled_prompt_embeds.repeat(1, num_images_per_prompt).view( | |
| bs_embed * num_images_per_prompt, -1 | |
| ) | |
| if do_classifier_free_guidance: | |
| negative_pooled_prompt_embeds = negative_pooled_prompt_embeds.repeat(1, num_images_per_prompt).view( | |
| bs_embed * num_images_per_prompt, -1 | |
| ) | |
| return prompt_embeds, negative_prompt_embeds, pooled_prompt_embeds, negative_pooled_prompt_embeds | |
| # Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion.StableDiffusionPipeline.prepare_extra_step_kwargs | |
| def prepare_extra_step_kwargs(self, generator, eta): | |
| # prepare extra kwargs for the scheduler step, since not all schedulers have the same signature | |
| # eta (η) is only used with the DDIMScheduler, it will be ignored for other schedulers. | |
| # eta corresponds to η in DDIM paper: https://arxiv.org/abs/2010.02502 | |
| # and should be between [0, 1] | |
| accepts_eta = "eta" in set(inspect.signature(self.scheduler.step).parameters.keys()) | |
| extra_step_kwargs = {} | |
| if accepts_eta: | |
| extra_step_kwargs["eta"] = eta | |
| # check if the scheduler accepts generator | |
| accepts_generator = "generator" in set(inspect.signature(self.scheduler.step).parameters.keys()) | |
| if accepts_generator: | |
| extra_step_kwargs["generator"] = generator | |
| return extra_step_kwargs | |
| def check_inputs( | |
| self, | |
| prompt, | |
| prompt_2, | |
| strength, | |
| num_inference_steps, | |
| callback_steps, | |
| negative_prompt=None, | |
| negative_prompt_2=None, | |
| prompt_embeds=None, | |
| negative_prompt_embeds=None, | |
| ): | |
| if strength < 0 or strength > 1: | |
| raise ValueError(f"The value of strength should in [0.0, 1.0] but is {strength}") | |
| if num_inference_steps is None: | |
| raise ValueError("`num_inference_steps` cannot be None.") | |
| elif not isinstance(num_inference_steps, int) or num_inference_steps <= 0: | |
| raise ValueError( | |
| f"`num_inference_steps` has to be a positive integer but is {num_inference_steps} of type" | |
| f" {type(num_inference_steps)}." | |
| ) | |
| if (callback_steps is None) or ( | |
| callback_steps is not None and (not isinstance(callback_steps, int) or callback_steps <= 0) | |
| ): | |
| raise ValueError( | |
| f"`callback_steps` has to be a positive integer but is {callback_steps} of type" | |
| f" {type(callback_steps)}." | |
| ) | |
| if prompt is not None and prompt_embeds is not None: | |
| raise ValueError( | |
| f"Cannot forward both `prompt`: {prompt} and `prompt_embeds`: {prompt_embeds}. Please make sure to" | |
| " only forward one of the two." | |
| ) | |
| elif prompt_2 is not None and prompt_embeds is not None: | |
| raise ValueError( | |
| f"Cannot forward both `prompt_2`: {prompt_2} and `prompt_embeds`: {prompt_embeds}. Please make sure to" | |
| " only forward one of the two." | |
| ) | |
| elif prompt is None and prompt_embeds is None: | |
| raise ValueError( | |
| "Provide either `prompt` or `prompt_embeds`. Cannot leave both `prompt` and `prompt_embeds` undefined." | |
| ) | |
| elif prompt is not None and (not isinstance(prompt, str) and not isinstance(prompt, list)): | |
| raise ValueError(f"`prompt` has to be of type `str` or `list` but is {type(prompt)}") | |
| elif prompt_2 is not None and (not isinstance(prompt_2, str) and not isinstance(prompt_2, list)): | |
| raise ValueError(f"`prompt_2` has to be of type `str` or `list` but is {type(prompt_2)}") | |
| if negative_prompt is not None and negative_prompt_embeds is not None: | |
| raise ValueError( | |
| f"Cannot forward both `negative_prompt`: {negative_prompt} and `negative_prompt_embeds`:" | |
| f" {negative_prompt_embeds}. Please make sure to only forward one of the two." | |
| ) | |
| elif negative_prompt_2 is not None and negative_prompt_embeds is not None: | |
| raise ValueError( | |
| f"Cannot forward both `negative_prompt_2`: {negative_prompt_2} and `negative_prompt_embeds`:" | |
| f" {negative_prompt_embeds}. Please make sure to only forward one of the two." | |
| ) | |
| if prompt_embeds is not None and negative_prompt_embeds is not None: | |
| if prompt_embeds.shape != negative_prompt_embeds.shape: | |
| raise ValueError( | |
| "`prompt_embeds` and `negative_prompt_embeds` must have the same shape when passed directly, but" | |
| f" got: `prompt_embeds` {prompt_embeds.shape} != `negative_prompt_embeds`" | |
| f" {negative_prompt_embeds.shape}." | |
| ) | |
| def get_timesteps(self, num_inference_steps, strength, device, denoising_start=None): | |
| # get the original timestep using init_timestep | |
| if denoising_start is None: | |
| init_timestep = min(int(num_inference_steps * strength), num_inference_steps) | |
| t_start = max(num_inference_steps - init_timestep, 0) | |
| else: | |
| t_start = 0 | |
| timesteps = self.scheduler.timesteps[t_start * self.scheduler.order :] | |
| # Strength is irrelevant if we directly request a timestep to start at; | |
| # that is, strength is determined by the denoising_start instead. | |
| if denoising_start is not None: | |
| discrete_timestep_cutoff = int( | |
| round( | |
| self.scheduler.config.num_train_timesteps | |
| - (denoising_start * self.scheduler.config.num_train_timesteps) | |
| ) | |
| ) | |
| timesteps = list(filter(lambda ts: ts < discrete_timestep_cutoff, timesteps)) | |
| return torch.tensor(timesteps), len(timesteps) | |
| return timesteps, num_inference_steps - t_start | |
| def prepare_latents( | |
| self, image, timestep, batch_size, num_images_per_prompt, dtype, device, generator=None, add_noise=True | |
| ): | |
| if not isinstance(image, (torch.Tensor, PIL.Image.Image, list)): | |
| raise ValueError( | |
| f"`image` has to be of type `torch.Tensor`, `PIL.Image.Image` or list but is {type(image)}" | |
| ) | |
| # Offload text encoder if `enable_model_cpu_offload` was enabled | |
| if hasattr(self, "final_offload_hook") and self.final_offload_hook is not None: | |
| self.text_encoder_2.to("cpu") | |
| torch.cuda.empty_cache() | |
| image = image.to(device=device, dtype=dtype) | |
| batch_size = batch_size * num_images_per_prompt | |
| if image.shape[1] == 4: | |
| init_latents = image | |
| else: | |
| # make sure the VAE is in float32 mode, as it overflows in float16 | |
| if self.vae.config.force_upcast: | |
| image = image.float() | |
| self.vae.to(dtype=torch.float32) | |
| if isinstance(generator, list) and len(generator) != batch_size: | |
| raise ValueError( | |
| f"You have passed a list of generators of length {len(generator)}, but requested an effective batch" | |
| f" size of {batch_size}. Make sure the batch size matches the length of the generators." | |
| ) | |
| elif isinstance(generator, list): | |
| init_latents = [ | |
| self.vae.encode(image[i : i + 1]).latent_dist.sample(generator[i]) for i in range(batch_size) | |
| ] | |
| init_latents = torch.cat(init_latents, dim=0) | |
| else: | |
| init_latents = self.vae.encode(image).latent_dist.sample(generator) | |
| if self.vae.config.force_upcast: | |
| self.vae.to(dtype) | |
| init_latents = init_latents.to(dtype) | |
| init_latents = self.vae.config.scaling_factor * init_latents | |
| if batch_size > init_latents.shape[0] and batch_size % init_latents.shape[0] == 0: | |
| # expand init_latents for batch_size | |
| additional_image_per_prompt = batch_size // init_latents.shape[0] | |
| init_latents = torch.cat([init_latents] * additional_image_per_prompt, dim=0) | |
| elif batch_size > init_latents.shape[0] and batch_size % init_latents.shape[0] != 0: | |
| raise ValueError( | |
| f"Cannot duplicate `image` of batch size {init_latents.shape[0]} to {batch_size} text prompts." | |
| ) | |
| else: | |
| init_latents = torch.cat([init_latents], dim=0) | |
| if add_noise: | |
| shape = init_latents.shape | |
| noise = randn_tensor(shape, generator=generator, device=device, dtype=dtype) | |
| # get latents | |
| init_latents = self.scheduler.add_noise(init_latents, noise, timestep) | |
| latents = init_latents | |
| return latents | |
| def _get_add_time_ids( | |
| self, original_size, crops_coords_top_left, target_size, aesthetic_score, negative_aesthetic_score, dtype | |
| ): | |
| if self.config.requires_aesthetics_score: | |
| add_time_ids = list(original_size + crops_coords_top_left + (aesthetic_score,)) | |
| add_neg_time_ids = list(original_size + crops_coords_top_left + (negative_aesthetic_score,)) | |
| else: | |
| add_time_ids = list(original_size + crops_coords_top_left + target_size) | |
| add_neg_time_ids = list(original_size + crops_coords_top_left + target_size) | |
| passed_add_embed_dim = ( | |
| self.unet.config.addition_time_embed_dim * len(add_time_ids) + self.text_encoder_2.config.projection_dim | |
| ) | |
| expected_add_embed_dim = self.unet.add_embedding.linear_1.in_features | |
| if ( | |
| expected_add_embed_dim > passed_add_embed_dim | |
| and (expected_add_embed_dim - passed_add_embed_dim) == self.unet.config.addition_time_embed_dim | |
| ): | |
| raise ValueError( | |
| f"Model expects an added time embedding vector of length {expected_add_embed_dim}, but a vector of {passed_add_embed_dim} was created. Please make sure to enable `requires_aesthetics_score` with `pipe.register_to_config(requires_aesthetics_score=True)` to make sure `aesthetic_score` {aesthetic_score} and `negative_aesthetic_score` {negative_aesthetic_score} is correctly used by the model." | |
| ) | |
| elif ( | |
| expected_add_embed_dim < passed_add_embed_dim | |
| and (passed_add_embed_dim - expected_add_embed_dim) == self.unet.config.addition_time_embed_dim | |
| ): | |
| raise ValueError( | |
| f"Model expects an added time embedding vector of length {expected_add_embed_dim}, but a vector of {passed_add_embed_dim} was created. Please make sure to disable `requires_aesthetics_score` with `pipe.register_to_config(requires_aesthetics_score=False)` to make sure `target_size` {target_size} is correctly used by the model." | |
| ) | |
| elif expected_add_embed_dim != passed_add_embed_dim: | |
| raise ValueError( | |
| f"Model expects an added time embedding vector of length {expected_add_embed_dim}, but a vector of {passed_add_embed_dim} was created. The model has an incorrect config. Please check `unet.config.time_embedding_type` and `text_encoder_2.config.projection_dim`." | |
| ) | |
| add_time_ids = torch.tensor([add_time_ids], dtype=dtype) | |
| add_neg_time_ids = torch.tensor([add_neg_time_ids], dtype=dtype) | |
| return add_time_ids, add_neg_time_ids | |
| # Copied from diffusers.pipelines.stable_diffusion.pipeline_stable_diffusion_upscale.StableDiffusionUpscalePipeline.upcast_vae | |
| def upcast_vae(self): | |
| dtype = self.vae.dtype | |
| self.vae.to(dtype=torch.float32) | |
| use_torch_2_0_or_xformers = isinstance( | |
| self.vae.decoder.mid_block.attentions[0].processor, | |
| ( | |
| AttnProcessor2_0, | |
| XFormersAttnProcessor, | |
| LoRAXFormersAttnProcessor, | |
| LoRAAttnProcessor2_0, | |
| ), | |
| ) | |
| # if xformers or torch_2_0 is used attention block does not need | |
| # to be in float32 which can save lots of memory | |
| if use_torch_2_0_or_xformers: | |
| self.vae.post_quant_conv.to(dtype) | |
| self.vae.decoder.conv_in.to(dtype) | |
| self.vae.decoder.mid_block.to(dtype) | |
| def __call__( | |
| self, | |
| prompt: Union[str, List[str]] = None, | |
| prompt_2: Optional[Union[str, List[str]]] = None, | |
| image: Union[ | |
| torch.FloatTensor, | |
| PIL.Image.Image, | |
| np.ndarray, | |
| List[torch.FloatTensor], | |
| List[PIL.Image.Image], | |
| List[np.ndarray], | |
| ] = None, | |
| strength: float = 0.3, | |
| num_inference_steps: int = 50, | |
| denoising_start: Optional[float] = None, | |
| denoising_end: Optional[float] = None, | |
| guidance_scale: float = 5.0, | |
| negative_prompt: Optional[Union[str, List[str]]] = None, | |
| negative_prompt_2: Optional[Union[str, List[str]]] = None, | |
| num_images_per_prompt: Optional[int] = 1, | |
| eta: float = 0.0, | |
| generator: Optional[Union[torch.Generator, List[torch.Generator]]] = None, | |
| latents: Optional[torch.FloatTensor] = None, | |
| prompt_embeds: Optional[torch.FloatTensor] = None, | |
| negative_prompt_embeds: Optional[torch.FloatTensor] = None, | |
| pooled_prompt_embeds: Optional[torch.FloatTensor] = None, | |
| negative_pooled_prompt_embeds: Optional[torch.FloatTensor] = None, | |
| output_type: Optional[str] = "pil", | |
| return_dict: bool = True, | |
| callback: Optional[Callable[[int, int, torch.FloatTensor], None]] = None, | |
| callback_steps: int = 1, | |
| cross_attention_kwargs: Optional[Dict[str, Any]] = None, | |
| guidance_rescale: float = 0.0, | |
| original_size: Tuple[int, int] = None, | |
| crops_coords_top_left: Tuple[int, int] = (0, 0), | |
| target_size: Tuple[int, int] = None, | |
| aesthetic_score: float = 6.0, | |
| negative_aesthetic_score: float = 2.5, | |
| map: torch.FloatTensor = None, | |
| original_image: Union[ | |
| torch.FloatTensor, | |
| PIL.Image.Image, | |
| np.ndarray, | |
| List[torch.FloatTensor], | |
| List[PIL.Image.Image], | |
| List[np.ndarray], | |
| ] = None, | |
| ): | |
| r""" | |
| Function invoked when calling the pipeline for generation. | |
| Args: | |
| prompt (`str` or `List[str]`, *optional*): | |
| The prompt or prompts to guide the image generation. If not defined, one has to pass `prompt_embeds`. | |
| instead. | |
| prompt_2 (`str` or `List[str]`, *optional*): | |
| The prompt or prompts to be sent to the `tokenizer_2` and `text_encoder_2`. If not defined, `prompt` is | |
| used in both text-encoders | |
| image (`torch.FloatTensor` or `PIL.Image.Image` or `np.ndarray` or `List[torch.FloatTensor]` or `List[PIL.Image.Image]` or `List[np.ndarray]`): | |
| The image(s) to modify with the pipeline. | |
| strength (`float`, *optional*, defaults to 0.3): | |
| Conceptually, indicates how much to transform the reference `image`. Must be between 0 and 1. `image` | |
| will be used as a starting point, adding more noise to it the larger the `strength`. The number of | |
| denoising steps depends on the amount of noise initially added. When `strength` is 1, added noise will | |
| be maximum and the denoising process will run for the full number of iterations specified in | |
| `num_inference_steps`. A value of 1, therefore, essentially ignores `image`. Note that in the case of | |
| `denoising_start` being declared as an integer, the value of `strength` will be ignored. | |
| num_inference_steps (`int`, *optional*, defaults to 50): | |
| The number of denoising steps. More denoising steps usually lead to a higher quality image at the | |
| expense of slower inference. | |
| denoising_start (`float`, *optional*): | |
| When specified, indicates the fraction (between 0.0 and 1.0) of the total denoising process to be | |
| bypassed before it is initiated. Consequently, the initial part of the denoising process is skipped and | |
| it is assumed that the passed `image` is a partly denoised image. Note that when this is specified, | |
| strength will be ignored. The `denoising_start` parameter is particularly beneficial when this pipeline | |
| is integrated into a "Mixture of Denoisers" multi-pipeline setup, as detailed in [**Refining the Image | |
| Output**](https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/stable_diffusion_xl#refining-the-image-output). | |
| denoising_end (`float`, *optional*): | |
| When specified, determines the fraction (between 0.0 and 1.0) of the total denoising process to be | |
| completed before it is intentionally prematurely terminated. As a result, the returned sample will | |
| still retain a substantial amount of noise (ca. final 20% of timesteps still needed) and should be | |
| denoised by a successor pipeline that has `denoising_start` set to 0.8 so that it only denoises the | |
| final 20% of the scheduler. The denoising_end parameter should ideally be utilized when this pipeline | |
| forms a part of a "Mixture of Denoisers" multi-pipeline setup, as elaborated in [**Refining the Image | |
| Output**](https://huggingface.co/docs/diffusers/api/pipelines/stable_diffusion/stable_diffusion_xl#refining-the-image-output). | |
| guidance_scale (`float`, *optional*, defaults to 7.5): | |
| Guidance scale as defined in [Classifier-Free Diffusion Guidance](https://arxiv.org/abs/2207.12598). | |
| `guidance_scale` is defined as `w` of equation 2. of [Imagen | |
| Paper](https://arxiv.org/pdf/2205.11487.pdf). Guidance scale is enabled by setting `guidance_scale > | |
| 1`. Higher guidance scale encourages to generate images that are closely linked to the text `prompt`, | |
| usually at the expense of lower image quality. | |
| negative_prompt (`str` or `List[str]`, *optional*): | |
| The prompt or prompts not to guide the image generation. If not defined, one has to pass | |
| `negative_prompt_embeds` instead. Ignored when not using guidance (i.e., ignored if `guidance_scale` is | |
| less than `1`). | |
| negative_prompt_2 (`str` or `List[str]`, *optional*): | |
| The prompt or prompts not to guide the image generation to be sent to `tokenizer_2` and | |
| `text_encoder_2`. If not defined, `negative_prompt` is used in both text-encoders | |
| num_images_per_prompt (`int`, *optional*, defaults to 1): | |
| The number of images to generate per prompt. | |
| eta (`float`, *optional*, defaults to 0.0): | |
| Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to | |
| [`schedulers.DDIMScheduler`], will be ignored for others. | |
| generator (`torch.Generator` or `List[torch.Generator]`, *optional*): | |
| One or a list of [torch generator(s)](https://pytorch.org/docs/stable/generated/torch.Generator.html) | |
| to make generation deterministic. | |
| latents (`torch.FloatTensor`, *optional*): | |
| Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image | |
| generation. Can be used to tweak the same generation with different prompts. If not provided, a latents | |
| tensor will ge generated by sampling using the supplied random `generator`. | |
| prompt_embeds (`torch.FloatTensor`, *optional*): | |
| Pre-generated text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. If not | |
| provided, text embeddings will be generated from `prompt` input argument. | |
| negative_prompt_embeds (`torch.FloatTensor`, *optional*): | |
| Pre-generated negative text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt | |
| weighting. If not provided, negative_prompt_embeds will be generated from `negative_prompt` input | |
| argument. | |
| pooled_prompt_embeds (`torch.FloatTensor`, *optional*): | |
| Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt weighting. | |
| If not provided, pooled text embeddings will be generated from `prompt` input argument. | |
| negative_pooled_prompt_embeds (`torch.FloatTensor`, *optional*): | |
| Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, *e.g.* prompt | |
| weighting. If not provided, pooled negative_prompt_embeds will be generated from `negative_prompt` | |
| input argument. | |
| output_type (`str`, *optional*, defaults to `"pil"`): | |
| The output format of the generate image. Choose between | |
| [PIL](https://pillow.readthedocs.io/en/stable/): `PIL.Image.Image` or `np.array`. | |
| return_dict (`bool`, *optional*, defaults to `True`): | |
| Whether or not to return a [`~pipelines.stable_diffusion.StableDiffusionXLPipelineOutput`] instead of a | |
| plain tuple. | |
| callback (`Callable`, *optional*): | |
| A function that will be called every `callback_steps` steps during inference. The function will be | |
| called with the following arguments: `callback(step: int, timestep: int, latents: torch.FloatTensor)`. | |
| callback_steps (`int`, *optional*, defaults to 1): | |
| The frequency at which the `callback` function will be called. If not specified, the callback will be | |
| called at every step. | |
| cross_attention_kwargs (`dict`, *optional*): | |
| A kwargs dictionary that if specified is passed along to the `AttentionProcessor` as defined under | |
| `self.processor` in | |
| [diffusers.cross_attention](https://github.com/huggingface/diffusers/blob/main/src/diffusers/models/cross_attention.py). | |
| guidance_rescale (`float`, *optional*, defaults to 0.7): | |
| Guidance rescale factor proposed by [Common Diffusion Noise Schedules and Sample Steps are | |
| Flawed](https://arxiv.org/pdf/2305.08891.pdf) `guidance_scale` is defined as `φ` in equation 16. of | |
| [Common Diffusion Noise Schedules and Sample Steps are Flawed](https://arxiv.org/pdf/2305.08891.pdf). | |
| Guidance rescale factor should fix overexposure when using zero terminal SNR. | |
| original_size (`Tuple[int]`, *optional*, defaults to (1024, 1024)): | |
| If `original_size` is not the same as `target_size` the image will appear to be down- or upsampled. | |
| `original_size` defaults to `(width, height)` if not specified. Part of SDXL's micro-conditioning as | |
| explained in section 2.2 of | |
| [https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). | |
| crops_coords_top_left (`Tuple[int]`, *optional*, defaults to (0, 0)): | |
| `crops_coords_top_left` can be used to generate an image that appears to be "cropped" from the position | |
| `crops_coords_top_left` downwards. Favorable, well-centered images are usually achieved by setting | |
| `crops_coords_top_left` to (0, 0). Part of SDXL's micro-conditioning as explained in section 2.2 of | |
| [https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). | |
| target_size (`Tuple[int]`, *optional*, defaults to (1024, 1024)): | |
| For most cases, `target_size` should be set to the desired height and width of the generated image. If | |
| not specified it will default to `(width, height)`. Part of SDXL's micro-conditioning as explained in | |
| section 2.2 of [https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). | |
| aesthetic_score (`float`, *optional*, defaults to 6.0): | |
| Used to simulate an aesthetic score of the generated image by influencing the positive text condition. | |
| Part of SDXL's micro-conditioning as explained in section 2.2 of | |
| [https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). | |
| negative_aesthetic_score (`float`, *optional*, defaults to 2.5): | |
| Part of SDXL's micro-conditioning as explained in section 2.2 of | |
| [https://huggingface.co/papers/2307.01952](https://huggingface.co/papers/2307.01952). Can be used to | |
| simulate an aesthetic score of the generated image by influencing the negative text condition. | |
| Examples: | |
| Returns: | |
| [`~pipelines.stable_diffusion.StableDiffusionXLPipelineOutput`] or `tuple`: | |
| [`~pipelines.stable_diffusion.StableDiffusionXLPipelineOutput`] if `return_dict` is True, otherwise a | |
| `tuple. When returning a tuple, the first element is a list with the generated images. | |
| """ | |
| # 1. Check inputs. Raise error if not correct | |
| self.check_inputs( | |
| prompt, | |
| prompt_2, | |
| strength, | |
| num_inference_steps, | |
| callback_steps, | |
| negative_prompt, | |
| negative_prompt_2, | |
| prompt_embeds, | |
| negative_prompt_embeds, | |
| ) | |
| # 2. Define call parameters | |
| if prompt is not None and isinstance(prompt, str): | |
| batch_size = 1 | |
| elif prompt is not None and isinstance(prompt, list): | |
| batch_size = len(prompt) | |
| else: | |
| batch_size = prompt_embeds.shape[0] | |
| device = self._execution_device | |
| # here `guidance_scale` is defined analog to the guidance weight `w` of equation (2) | |
| # of the Imagen paper: https://arxiv.org/pdf/2205.11487.pdf . `guidance_scale = 1` | |
| # corresponds to doing no classifier free guidance. | |
| do_classifier_free_guidance = guidance_scale > 1.0 | |
| # 3. Encode input prompt | |
| text_encoder_lora_scale = ( | |
| cross_attention_kwargs.get("scale", None) if cross_attention_kwargs is not None else None | |
| ) | |
| ( | |
| prompt_embeds, | |
| negative_prompt_embeds, | |
| pooled_prompt_embeds, | |
| negative_pooled_prompt_embeds, | |
| ) = self.encode_prompt( | |
| prompt=prompt, | |
| prompt_2=prompt_2, | |
| device=device, | |
| num_images_per_prompt=num_images_per_prompt, | |
| do_classifier_free_guidance=do_classifier_free_guidance, | |
| negative_prompt=negative_prompt, | |
| negative_prompt_2=negative_prompt_2, | |
| prompt_embeds=prompt_embeds, | |
| negative_prompt_embeds=negative_prompt_embeds, | |
| pooled_prompt_embeds=pooled_prompt_embeds, | |
| negative_pooled_prompt_embeds=negative_pooled_prompt_embeds, | |
| lora_scale=text_encoder_lora_scale, | |
| ) | |
| # 4. Preprocess image | |
| #image = self.image_processor.preprocess(image) #ideally we would have preprocess the image with diffusers, but for this POC we won't --- it throws a deprecated warning | |
| map = torchvision.transforms.Resize(tuple(s // self.vae_scale_factor for s in original_image.shape[2:]),antialias=None)(map) | |
| # 5. Prepare timesteps | |
| def denoising_value_valid(dnv): | |
| return type(denoising_end) == float and 0 < dnv < 1 | |
| self.scheduler.set_timesteps(num_inference_steps, device=device) | |
| #begin diff diff change | |
| total_time_steps = num_inference_steps | |
| #end diff diff change | |
| timesteps, num_inference_steps = self.get_timesteps( | |
| num_inference_steps, strength, device, denoising_start=denoising_start if denoising_value_valid else None | |
| ) | |
| latent_timestep = timesteps[:1].repeat(batch_size * num_images_per_prompt) | |
| add_noise = True if denoising_start is None else False | |
| # 6. Prepare latent variables | |
| latents = self.prepare_latents( | |
| image, | |
| latent_timestep, | |
| batch_size, | |
| num_images_per_prompt, | |
| prompt_embeds.dtype, | |
| device, | |
| generator, | |
| add_noise, | |
| ) | |
| # 7. Prepare extra step kwargs. | |
| extra_step_kwargs = self.prepare_extra_step_kwargs(generator, eta) | |
| height, width = latents.shape[-2:] | |
| height = height * self.vae_scale_factor | |
| width = width * self.vae_scale_factor | |
| original_size = original_size or (height, width) | |
| target_size = target_size or (height, width) | |
| # 8. Prepare added time ids & embeddings | |
| add_text_embeds = pooled_prompt_embeds | |
| add_time_ids, add_neg_time_ids = self._get_add_time_ids( | |
| original_size, | |
| crops_coords_top_left, | |
| target_size, | |
| aesthetic_score, | |
| negative_aesthetic_score, | |
| dtype=prompt_embeds.dtype, | |
| ) | |
| add_time_ids = add_time_ids.repeat(batch_size * num_images_per_prompt, 1) | |
| if do_classifier_free_guidance: | |
| prompt_embeds = torch.cat([negative_prompt_embeds, prompt_embeds], dim=0) | |
| add_text_embeds = torch.cat([negative_pooled_prompt_embeds, add_text_embeds], dim=0) | |
| add_neg_time_ids = add_neg_time_ids.repeat(batch_size * num_images_per_prompt, 1) | |
| add_time_ids = torch.cat([add_neg_time_ids, add_time_ids], dim=0) | |
| prompt_embeds = prompt_embeds.to(device) | |
| add_text_embeds = add_text_embeds.to(device) | |
| add_time_ids = add_time_ids.to(device) | |
| # 9. Denoising loop | |
| num_warmup_steps = max(len(timesteps) - num_inference_steps * self.scheduler.order, 0) | |
| # 9.1 Apply denoising_end | |
| if ( | |
| denoising_end is not None | |
| and denoising_start is not None | |
| and denoising_value_valid(denoising_end) | |
| and denoising_value_valid(denoising_start) | |
| and denoising_start >= denoising_end | |
| ): | |
| raise ValueError( | |
| f"`denoising_start`: {denoising_start} cannot be larger than or equal to `denoising_end`: " | |
| + f" {denoising_end} when using type float." | |
| ) | |
| elif denoising_end is not None and denoising_value_valid(denoising_end): | |
| discrete_timestep_cutoff = int( | |
| round( | |
| self.scheduler.config.num_train_timesteps | |
| - (denoising_end * self.scheduler.config.num_train_timesteps) | |
| ) | |
| ) | |
| num_inference_steps = len(list(filter(lambda ts: ts >= discrete_timestep_cutoff, timesteps))) | |
| timesteps = timesteps[:num_inference_steps] | |
| # prepartions for diff diff | |
| original_with_noise = self.prepare_latents( | |
| original_image, timesteps, batch_size, num_images_per_prompt, prompt_embeds.dtype, device, generator | |
| ) | |
| thresholds = torch.arange(total_time_steps, dtype=map.dtype) / total_time_steps | |
| thresholds = thresholds.unsqueeze(1).unsqueeze(1).to(device) | |
| masks = map > (thresholds + (denoising_start or 0)) | |
| # end diff diff preparations | |
| with self.progress_bar(total=num_inference_steps) as progress_bar: | |
| for i, t in enumerate(timesteps): | |
| # diff diff | |
| if i==0 and denoising_start is None: | |
| latents = original_with_noise[:1] | |
| else: | |
| mask = masks[i].unsqueeze(0) | |
| # cast mask to the same type as latents etc | |
| mask = mask.to(latents.dtype) | |
| mask = mask.unsqueeze(1) # fit shape | |
| latents = original_with_noise[i] * mask + latents * (1 - mask) | |
| # end diff diff | |
| # expand the latents if we are doing classifier free guidance | |
| latent_model_input = torch.cat([latents] * 2) if do_classifier_free_guidance else latents | |
| latent_model_input = self.scheduler.scale_model_input(latent_model_input, t) | |
| # predict the noise residual | |
| added_cond_kwargs = {"text_embeds": add_text_embeds, "time_ids": add_time_ids} | |
| noise_pred = self.unet( | |
| latent_model_input, | |
| t, | |
| encoder_hidden_states=prompt_embeds, | |
| cross_attention_kwargs=cross_attention_kwargs, | |
| added_cond_kwargs=added_cond_kwargs, | |
| return_dict=False, | |
| )[0] | |
| # perform guidance | |
| if do_classifier_free_guidance: | |
| noise_pred_uncond, noise_pred_text = noise_pred.chunk(2) | |
| noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond) | |
| if do_classifier_free_guidance and guidance_rescale > 0.0: | |
| # Based on 3.4. in https://arxiv.org/pdf/2305.08891.pdf | |
| noise_pred = rescale_noise_cfg(noise_pred, noise_pred_text, guidance_rescale=guidance_rescale) | |
| # compute the previous noisy sample x_t -> x_t-1 | |
| latents = self.scheduler.step(noise_pred, t, latents, **extra_step_kwargs, return_dict=False)[0] | |
| # call the callback, if provided | |
| if i == len(timesteps) - 1 or ((i + 1) > num_warmup_steps and (i + 1) % self.scheduler.order == 0): | |
| progress_bar.update() | |
| if callback is not None and i % callback_steps == 0: | |
| callback(i, t, latents) | |
| # make sure the VAE is in float32 mode, as it overflows in float16 | |
| if self.vae.dtype == torch.float16 and self.vae.config.force_upcast: | |
| self.upcast_vae() | |
| latents = latents.to(next(iter(self.vae.post_quant_conv.parameters())).dtype) | |
| if not output_type == "latent": | |
| image = self.vae.decode(latents / self.vae.config.scaling_factor, return_dict=False)[0] | |
| else: | |
| image = latents | |
| return StableDiffusionXLPipelineOutput(images=image) | |
| # apply watermark if available | |
| if self.watermark is not None: | |
| image = self.watermark.apply_watermark(image) | |
| image = self.image_processor.postprocess(image, output_type=output_type) | |
| # Offload last model to CPU | |
| if hasattr(self, "final_offload_hook") and self.final_offload_hook is not None: | |
| self.final_offload_hook.offload() | |
| if not return_dict: | |
| return (image,) | |
| return StableDiffusionXLPipelineOutput(images=image) |